Stable diffusion base model. Mar 24, 2023 · December 7, 2022.

safetensors, it needed to use relative paths (Checkpoints\Checkpoints\01 - Photorealistic\model Use this model. The main work of the Base model is consistent with that of Stable Diffusion, with the ability to perform text-to-image, image-to-image, and image inpainting. History:13 commits. New stable diffusion model ( Stable Diffusion 2. Feb 22, 2023 · L'entraînement supplémentaire est réalisé en entraînant un modèle de base avec un ensemble de données supplémentaire qui vous intéresse. Stable Diffusion is a text-to-image model that generates photo-realistic images given any text input. For SDXL : Same as above. 5 or SD v2. Figure 1: Imagining mycelium couture. Jul 31, 2023 · Check out the Quick Start Guide if you are new to Stable Diffusion. Fine-Tuned Models (ie any checkpoint you download from CivitAI) = College. Stable Diffusion. Following the successful release of Stable Diffusion XL beta in April, SDXL 0. 0, generates high quality photorealsitic images, offers vibrant, accurate colors, superior contrast, and detailed shadows than the base SDXL at a native resolution of 1024x1024. 4 contributors. 0. 1 to use it with 2. TensorRT INT8 quantization is available now, with FP8 expected soon. 1-v, Hugging Face) at 768x768 resolution and (Stable Diffusion 2. This weights here are intended to be used with the 🧨 Nov 26, 2022 · Hi there :) I need to move Models directory to a separate 2TB drive to create some space on the iMac so followed these instructions for command line args. This model is perfect for generating anime-style images of characters, objects, animals, landscapes, and more. Jul 27, 2023 · This is interesting: I previously had my SDXL models (base + refiner) stored inside a subdirectory named "SDXL" under /models/Stable-Diffusion. These custom models usually perform better than the base models. Nov 21, 2023 · Using the Pick-a-Pic dataset of 851K crowdsourced pairwise preferences, we fine-tune the base model of the state-of-the-art Stable Diffusion XL (SDXL)-1. Model Name: Base Model | Model ID: base-model | Plug and play API's to generate images with Base Model. Jul 9, 2023 · Last update 10-01-2023 本記事について 概要 独自の基準で選んだ、Stable Diffusion v2モデル（と、TI embeddings）を紹介します。 下記の記事もお役に立てたら幸いです。 → Stable Diffusion v1モデル_H2-2023  → Stable Diffusion XLモデル_H2-2023  → ChatGPTでStable Diffusion用のプロンプトを作成する  SDv2 Jun 22, 2023 · This gives rise to the Stable Diffusion architecture. If you're not into that, then no. HassanBlend 1. Dec 24, 2023 · Stable Diffusion XL consisting of a Base model and a Refiner model. The Stable-Diffusion-Inpainting was initialized with the weights of the Stable-Diffusion-v-1-2. We are releasing Stable Video Diffusion, an image-to-video model, for research purposes: SVD: This model was trained to generate 14 frames at resolution 576x1024 given a context frame of the same size. 9, the most advanced development in the Stable Diffusion text-to-image suite of models. This Version is better suited for realism but also it handles drawings better. Stable Diffusion Inpainting is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input, with the extra capability of inpainting the pictures by using a mask. like 10. It is one of the best open-source weights provided by OpenCLIP. 9: The base model was trained on a variety of aspect ratios on images with resolution 1024^2. Optimum Optimum provides a Stable Diffusion pipeline compatible with both OpenVINO and ONNX Runtime . 31. Stable Diffusion 2 is based on OpenCLIP-ViT/H as the text-encoder, while the older architecture uses OpenAI’s ViT-L/14. Realistic Vision is the best Stable Diffusion model for generating realistic humans. For anime images, it is common to adjust Clip Skip and VAE settings based on the model you use. 4, released by Stability. Dec 16, 2023 · Prompt: “a black and white photo of a woman wearing a floral crown and holding a bouquet of flowers in the style of Bella Kotak” (base Stable Diffusion model) After fine-tuning the base model. The model is designed to be perfect for drawing using natural language descriptions. For a full list of model_id values and which models are fine-tunable, refer to Built-in Algorithms with pre-trained Model Table . Running on CPU Upgrade Apr 24, 2024 · LandscapeSuperMix. Highly accessible: It runs on a consumer grade Apr 28, 2024 · Base Model: SD 1. Apr 26, 2023 · A few months ago we showed how the MosaicML platform makes it simple—and cheap—to train a large-scale diffusion model from scratch. But people keeping 2. The model can be accessed via ClipDrop today, with API Aug 10, 2023 · 下載好後把 Base 跟 Refiner 丟到 \stable-diffusion-webui\models\Stable-diffusion 下面，VAE 丟到 \stable-diffusion-webui\models\VAE 下面。 丟好後等 WebUI 第一次跑完，連終端機整個關掉再重開一次即可，可以看到 checkpoint 有剛剛載好的模型。 What Can You Do with the Base Stable Diffusion Model? The base models of Stable Diffusion, such as XL 1. 9 produces massively improved image and composition detail over its predecessor. 0, on a less restrictive NSFW filtering of the LAION-5B dataset. ckpt) with 220k extra steps taken, with punsafe=0. Resources for more information: GitHub Jan 11, 2024 · Checkpoints like Copax Timeless SDXL, Zavychroma SDXL, Dreamshaper SDXL, Realvis SDXL, Samaritan 3D XL are fine-tuned on base SDXL 1. Explore thousands of high-quality Stable Diffusion models, share your AI-generated art, and engage with a vibrant community of creators For 1. (SVD) Image-to-Video is a latent diffusion model trained to generate short video clips from an image conditioning. On the Settings page, click User Interface on the left panel. Use it with the stablediffusion repository: download the v2-1_512-ema-pruned. You have your general-purpose liberal arts majors like Deliberate, Dreamshaper, or Lyriel. It’s so good at generating faces and eyes that it’s often hard to tell if the image is AI-generated. 4/1. stable-diffusion. Choose from thousands of models like Base Model or upload your custom models for free Base Model | Stable Diffusion API - Generate Unlimited Images Using Base Model Configuration: Stable Diffusion XL 1. 0 Reprinted models are for communication and learning purposes only, not for Feb 8, 2024 · Stable Diffusion Web UIで「モデル」を変更する方法を解説しています。「Civitai」などのサイトで公開されているモデルファイルをダウンロードして所定のフォルダに格納するだけで、簡単にモデルを変更できます。 Apr 11, 2024 · The XSArchi_127新科幻Neo Sci-Fi model on Civitai is a Stable Diffusion LoRA model that contains all sci-fi scenarios and subdivision styles. The model is an ensemble of experts pipeline, where the base model generates latents that are then further Hey guys, I'm training a lot of models lately using sd1. Mr-Jay on Oct 4, 2022. ckpt) with an additional 55k steps on the same dataset (with punsafe=0. Oct 5, 2022 · Go to last tab "Settings" and there on the bottom you will have option to chose your model as on pic example: After selecting make sure to Apply Settings and then restart whole program. Its based on my new not yet published DEMONCORE V4 "NeoDEMON". Model Description: This is a model that can be used to generate and modify images based on text prompts. Load SDXL refiner 1. Stable Diffusion consists of three parts: A text encoder, which turns your prompt into a latent vector. Thankfully by fine-tuning the base Stable Diffusion model using captioned images, the ability of the base model to generate better-looking pictures The Stable-Diffusion-v1-4 checkpoint was initialized with the weights of the Stable-Diffusion-v1-2 checkpoint and subsequently fine-tuned on 225k steps at resolution 512x512 on "laion-aesthetics v2 5+" and 10% dropping of the text-conditioning to improve classifier-free guidance sampling. The LandscapeSuperMix model, with the version number v2. 5 checkpoint = High School. First of all you want to select your Stable Diffusion checkpoint, also known as a model. Jun 23, 2024 · Version 10B "NeoDEMON" (Experimental Trained) This version is a complete rebuild based on the dataset of Version 5. 0 as our base teacher model and have trained on the LAION Art Aesthetic dataset with image scores above 7. Stable Diffusion XL (SDXL) is a powerful text-to-image generation model that iterates on the previous Stable Diffusion models in three key ways: the UNet is 3x larger and SDXL combines a second text encoder (OpenCLIP ViT-bigG/14) with the original text encoder to significantly increase the number of parameters. These models, designed to convert text prompts into images, offer general-p Mar 18, 2024 · November 21, 2023. The generative artificial intelligence technology is the premier product of Stability AI and is considered to be a part of the ongoing artificial intelligence boom . Jun 22, 2023 · 22 Jun. It is a Latent Diffusion Model that uses two fixed, pretrained text encoders (OpenCLIP-ViT/G and CLIP-ViT/L). Uploaded the best CKPT IMO at 88210 Steps. main. Then I started reading tips and tricks, joined several Discord servers, and then went full hands-on to train and fine-tuning my own models. The model is updated quite regularly and so many improvements have been made since its launch. Realistic Vision. 0/2. Oct 7, 2023 · 2. Step 2. Feb 22, 2024 · The Stable Diffusion 3 suite of models currently ranges from 800M to 8B parameters. Today, we are excited to show the results of our own training run: under $50k to train Stable Diffusion 2 base1 from scratch in 7. Improve the Results with Refiner. Once you graduate, there's little reason to go back. 5 as a base on different people and objects for photorealism. EpiCPhotoGasm: The Photorealism Prodigy. The Swift package relies on the Core ML model files generated by python_coreml_stable_diffusion. 探讨AI绘画中大模型对作品影响的文章，提供模型下载与说明。 Mar 8, 2024 · Below find a quick summary of the Best Stable Diffusion Models. Upload a set of images depicting a person, animal, object or art style you want to imitate. 1 Lora. The Base model consists of three modules: U-Net, VAE, and two CLIP Text Encoders. 1. ai, the creators of Stable Diffusion, in August 2022. . 4 and v1. Started with the basics, running the base model on HuggingFace, testing different prompts. 3. patrickvonplaten. ckpt here. Full fine-tuning of larger models (consisting of billions of parameters) is inherently expensive and time-consuming. 0 model along with its refiner model. Unlike the paper, we have chosen to train the two models on 1M images for 100K steps for the Small and 125K steps for the Tiny mode respectively. It's good for creating fantasy, anime and semi-realistic images. Using the prompt. Nov 6, 2023 · The first public base model was SD v1. The text-conditional model is then trained in the. This repository comprises: StableDiffusion, a Swift package that developers can add to their Xcode projects as a dependency to deploy image generation capabilities in their apps. If you run into issues during installation or runtime, please refer to Learn how to access the Stable Diffusion model online and locally by following the How to Run Stable Diffusion tutorial. You can construct an image generation workflow by chaining different blocks (called nodes) together. Stable Cascade achieves a compression factor of 42, meaning that it is possible to encode a. 4 (Photorealism) + Protogen x5. 1-base, HuggingFace) at 512x512 resolution, both based on the same number of parameters and architecture as 2. Choose. At FP16 precision, the size of the Mar 19, 2024 · Stable Diffusion Models, or checkpoint models, are pre-trained Stable Diffusion weights for generating a particular style of images. safetensors. It contains all the baseline knowledge for how to turn text into images. Browse base model Stable Diffusion models, checkpoints, hypernetworks, textual inversions, embeddings, Aesthetic Gradients, and LORAs Mar 4, 2024 · The array of fine-tuned Stable Diffusion models is abundant and ever-growing. CHECK "ABOUT THIS VERSION" ON THE RIGHT IF YOU ARE NOT ON "V6" FOR IMPORTANT INFORMATION. This model was trained to generate 25 frames at resolution 576x1024 given a context frame of the same size, finetuned from SVD Image-to-Video [14 frames] . The stable-diffusion-xl-base-1. Robin Rombach add weights. Fix deprecated float16/fp16 variant loading through new `version` API. Aug 28, 2023 · Today, most custom models are built on top of either SD v1. Stable Diffusion is a latent diffusion model conditioned on the (non-pooled) text embeddings of a CLIP ViT-L/14 text encoder. AI美女を生成するのにおすすめのモデルを紹介します。 こちらで紹介するのは日本人(アジア人)の美女に対応しているモデルですが、もし日本人っぽくならない場合は「Japanese actress」「Korean idol」といったプロンプトを入れるのがおすすめです。 Mar 18, 2024 · We are releasing two new diffusion models for research purposes: SDXL-base-0. 0 base on 1,126,400,000 images at 256x256 resolution and 1,740,800,000 images at 512x512 resolution. Protogen x3. 0 and the larger SDXL-1. 1 ), and then fine-tuned for another 155k extra steps with punsafe=0. Art & Eros (aEros) + RealEldenApocalypse by aine_captain. Model type: Diffusion-based text-to-image generative model. 5; Download Count: 504K; Reviews: Overwhelmingly Positive (100%) File Size: 1. 5 avec un ensemble de données supplémentaire de voitures vintage pour biaiser l'esthétique des voitures vers le sous-genre. This is Part 2 of the Stable Diffusion for Beginner's series. Ema是一种算法，可以近似获得近n步权重的平均值。. To aid your selection, we present a list of versatile models, from the widely celebrated Stable diffusion v1. highly compressed latent space. Generate the image with the base SDXL model. This model card focuses on the model associated with the Stable Diffusion v2 model, available here. ViT/H is trained on LAION-2B with an accuracy of 78. Jan 17, 2023 · Handpicked image at my liking (trying to cover all possible angle and expression) Each Person has 45 Image equal, so thats 9*45 Images. This stable-diffusion-2 model is resumed from stable-diffusion-2-base ( 512-base-ema. i'll try to upload a lower steps CKPT or higher. The train_dreambooth_lora_sdxl. 4. 0 using diffusion pipeline. Stable Diffusion v2 Model Card. EpiCPhotoGasm. [ [open-in-colab]] Stable Diffusion XL (SDXL) is a powerful text-to-image generation model that iterates on the previous Stable Diffusion models in three key ways: the UNet is 3x larger and SDXL combines a second text encoder (OpenCLIP ViT-bigG/14) with the original text encoder to significantly increase the number of May 16, 2024 · Once you’ve uploaded your image to the img2img tab we need to select a checkpoint and make a few changes to the settings. py script shows how to implement the training procedure and adapt it for Stable Diffusion XL. 98 on the same dataset. Here I will be using the revAnimated model. 0 / sd_xl_base_1. May 13, 2024 · Pony Diffusion V6 is a versatile SDXL finetune capable of producing stunning SFW and NSFW visuals of various anthro, feral, or humanoids species and their interactions based on simple natural language prompts. 5, to base inpainting model you get new impainting model that inpaints with this other model concepts trained. May 28, 2024 · Model overview. 0 by sviasem. using 100 Epoch. LoRA works by adding a smaller number of new weights to the stable-diffusion-2-1 模型是根据 stable-diffusion-2 (768-v-ema. This approach aims to align with our core values and democratize access, providing users with a variety of options for scalability and quality to best meet their creative needs. It works nicely for realistic faces and imitates authentic camera photos. It usually takes just a few minutes. to get started. 1024x1024 image to 24x24, while maintaining crisp reconstructions. 98. Our time estimates are based on training Stable Diffusion 2. Par exemple, vous pouvez entraîner Stable Diffusion v1. The base model uses OpenCLIP-ViT/G and CLIP-ViT/L for text encoding whereas the refiner model only uses the OpenCLIP model. It allows users to invoke the desired style data by typing a prompt. DreamBooth is a method to personalize text2image models like stable diffusion given just a few (3~5) images of a subject. We use the standard image encoder from SD 2. 0 model is a text-to-image generative AI model developed by Stability AI. 5, because of their high quality image descriptions. ComfyUI breaks down a workflow into rearrangeable elements so you can easily make your own. Mar 24, 2023 · December 7, 2022. 💡 Note: For now, we only allow DreamBooth fine-tuning of the SDXL UNet via LoRA. Stable Diffusion XL. You’ll also find suggested models that will help you create the various styles on this list. Feb 12, 2024 · 2. It is a landscape-focused model that can generate various types of landscapes, including urban, architectural, and natural scenes. Sep 14, 2023 · Stable Diffusion XL(SDXL)とは、Stability AI 社が開発した最新のAI画像生成モデルです。以前のモデルに比べて、細かい部分もしっかりと反映してくれるようになり、より高画質なイラストを生成してくれます。そんなSDXLの導入方法・使い方について解説しています。 No code required to produce your model! Step 1. 0 model with Diffusion-DPO. I’ve been playing around with Stable Diffusion for some weeks now. Prompt: oil painting of zwx in style of van gogh. 3 (Photorealism) by darkstorm2150. Feb 20, 2023 · The following code shows how to fine-tune a Stable Diffusion 2. encoded to 128x128. For all the example images shared below, I’ve used the Stable Diffusion XL 1. 1 base model identified by model_id model-txt2img-stabilityai-stable-diffusion-v2-1-base on a custom training dataset. 5. Version 2. Juggernaut XL: Overall best Stable Diffusion model. SD 1. 2 by sdhassan. This model card focuses on the model associated with the Stable Diffusion v2-1-base model. To further improve the image quality and model accuracy, we will use Refiner. No token limit for prompts (original stable diffusion lets you use up to 75 tokens) DeepDanbooru integration, creates danbooru style tags for anime prompts xformers , major speed increase for select cards: (add --xformers to commandline args) This model card focuses on the model associated with the Stable Diffusion v2-1 model, codebase available here. 5: To use it as training base. 0 or the newer SD 3. 45 days using the MosaicML platform. Resumed for another 140k steps on 768x768 images. ckpt) and trained for 150k steps using a v-objective on the same dataset. 5k. stable-diffusion-2-1-base. 1 or v1. A decoder, which turns the final 64x64 latent patch into a higher-resolution 512x512 image. Realism Engine SDXL: Best Stable Diffusion model for photorealism. This stable-diffusion-2-1-base model fine-tunes stable-diffusion-2-base (512-base-ema. 5 model, but luckily by adding weight difference between other model and 1. Install the Models: Find the installation directory of the software you’re using to work with stable diffusion models. safetensors). 98。. 上面2个版本分别为包含 EMA 与不包含 EMA 。. Initially there was only one inpainting model - trained for base 1. ckpt) 进行微调的，在同一数据集上增加了 55k 个步骤（punsafe=0. 1: Same as above. co/stabilityai/stable-diffusion-xl-base-1. What kind of images a model generates depends on the training images. Parts of the graphics are from my Hephaistos 3. 0, and SD v2. Trained on dreambooth local without EMA using nonpruned 1,5 SD base model, fp16, Xformers. 5, SD v2. 1. A model won’t be able to generate a cat’s image if there’s never a cat in the training data. New stable diffusion model (Stable Diffusion 2. Aug 1, 2023 · We have taken Realistic-Vision 4. Analog Sep 25, 2023 · Stable Diffusionの実写・リアル系おすすめモデル. Nov 29, 2022 · Text encoder. oil painting of zwx in style of van gogh. If the model is in a subfolder, like I was using: C:\AI\stable-diffusion-webui\models\Stable-diffusion\Checkpoints\Checkpoints\01 - Photorealistic\model. Counterfeit is one of the most popular anime models for Stable Diffusion and has over 200K downloads. 1, but replace the decoder with a temporally-aware deflickering decoder. Our cost estimates are based on $2 / A100-hour. Jul 26, 2023 · Use this model main stable-diffusion-xl-base-1. Uber Realistic Porn Merge (URPM) by saftle. 0 base model; images resolution=1024×1024; Batch size=1; Euler scheduler for 50 steps; NVIDIA RTX 6000 Ada GPU. If the model is in the checkpoint directory, it just needs the model-name (model. 1, is a stable diffusion checkpoint available on Civitai. Our fine-tuned base model significantly outperforms both base SDXL-1. (it's also better to use other finetuned model as base) For 2. While I already like the… Jul 27, 2023 · Model reprinted from : https://huggingface. Model. In the Quicksetting List, add the following. A diffusion model, which repeatedly "denoises" a 64x64 latent image patch. history stable-diffusion-inpainting. The model and the code that uses the model to generate the image (also known as inference code). Jan 6, 2023 · ※複数のv1系Stable Diffusionモデルにおいて、データの破損（特定の位置にあるトークンが無視される現象）が確認されています。 それらには目印「 」を付けておきますので、下記の記事で詳細を確認してください。 Best Stable Diffusion Models - PhotoRealistic Styles. DreamShaper XL: Best alternative to Midjourney. ( #13) 5ede9e4 about 1 year ago. Jan 4, 2024 · Stable Diffusion model has been extensively employed in the study of architectural image generation, but there is still an opportunity to enhance in terms of the controllability of the generated 知乎专栏提供一个平台，让用户可以随心所欲地写作和自由表达自己的观点。 2 days ago · LoRA, or Low-Rank Adaptation, is a lightweight training technique used for fine-tuning Large Language and Stable Diffusion Models without needing full model training. You can also combine it with LORA models to be more versatile and generate unique artwork. with my newly trained model, I am happy with what I got: Images from dreambooth model. May 12, 2024 · Without them it would not have been possible to create this model. 5 models, each with their unique allure and general-purpose capabilities, to the SDXL model, a veritable upgrade boasting higher resolutions and quality. 1），然后针对 另外 155k 额外步骤，punsafe=0. Apr 16, 2023 · There are two primary techniques for fine-tuning, namely, additional training and using Dreambooth extension, both of which begin with a base model such as Stable Diffusion v2. 3. LandscapeSuperMix is a Stable Diffusion checkpoint model for cityscape. 0, are versatile tools capable of generating a broad spectrum of images across various styles, from photorealistic to animated and digital art. In the following weeks and months, they released SD v1. What It Does: Highly tuned for photorealism, this model excels in creating realistic images with minimal prompting. Jan 27, 2024 · For every Stable Diffusion style shared below, I’ve listed the prompts and two example images I’ve generated. 0 model consisting of an additional refinement model in human evaluation Feb 15, 2024 · Stable Diffusion uses a compression factor of 8, resulting in a 1024x1024 image being. You can even make your May 5, 2023 · Ecotech City, by Stable Diffusion. This stable-diffusion-2-1 model is fine-tuned from stable-diffusion-2 ( 768-v-ema. Today, Stability AI announces SDXL 0. Dreamlike Photoreal 2. Feb 15, 2024 · For more information on how to use Stable Diffusion XL with diffusers, please have a look at the Stable Diffusion XL Docs. "Stable Diffusion model" is used to refer to the official base models by StabilityAI, as well as all of these custom models. Train. Stable Diffusion 3 combines a diffusion transformer architecture and flow matching. It is a Latent Diffusion Model that uses two fixed, pretrained text encoders ( OpenCLIP-ViT/G and CLIP-ViT/L ). What makes Stable Diffusion unique ? It is completely open source. It handles various ethnicities and ages with ease. Copy the Model Files: Copy the downloaded model files from the downloads directory and paste them into the “models” directory of the software. f298da3 12 months ago. 0 and fine-tuned on 2. Jul 6, 2024 · ComfyUI is a node-based GUI for Stable Diffusion. Some commonly used blocks are Loading a Checkpoint Model, entering a prompt, specifying a sampler, etc. Wait for the custom stable diffusion model to be trained. Model Description. We provide a reference script for sampling, but there also exists a diffusers integration, which we expect to see more active community development. HASDX by bestjammer. We also finetune the widely used f8-decoder for temporal Sep 15, 2023 · Developed by: Stability AI. It has been trained for 4,000 epochs. The benchmark for TensorRT FP8 may change upon release. Most people will keep it until some better finetuned models show up. Overview. 1-v, Hugging Face) at 768x768 resolution and ( Stable Diffusion 2. Dec 11, 2023 · Stable Diffusionのモデルの中でも世界最高峰の画像生成モデルとして利用者から高い評価を受けている「SDXL」対応モデル。そんな話題のSDXL対応モデルのおすすめを人気順に20選まとめています！是非この記事をチェックして、最新モデルのクオリティの高さを実感してみてください。 Mar 10, 2024 · Apr 29, 2023. Nov 12, 2023 · Stable Diffusion is an artificial intelligence (AI) model that is revolutionising the world of digital art, making it easy for artists and content creators to generate high-quality images simply by using text prompts and source images. It is convenient to enable them in Quick Settings. Although stable diffusion released the base model, there have been many more pruned models released in recent months, and other models such as a lora and embeddings Creating an Image To create an image using Stable Diffusion, you'll typically follow a process involving setting up the necessary software environment, obtaining the model, and then We benchmarked the U-Net training throughput as we scale the number of A100 GPUs from 8 to 128. AniVerse: Best Stable Diffusion model for anime. download Copy download link. 该 Although stable diffusion released the base model, there have been many more pruned models released in recent months, and other models such as a lora and embeddings Creating an Image To create an image using Stable Diffusion, you'll typically follow a process involving setting up the necessary software environment, obtaining the model, and then Stable Diffusion is a deep learning, text-to-image model released in 2022 based on diffusion techniques. Text-to-Image with Stable Diffusion. 1 reply. 99GB SafeTensor; License: CreativeML Open RAIL-M Addendum; EpiCRealism is another top-tier Stable Diffusion model for photorealism generation. 2. Jan 17, 2024 · Step 4: Testing the model (optional) You can also use the second cell of the notebook to test using the model. AI generators utilise this model to give you photorealistic images and other detailed digital illustrations at the click of a button. if ik ni gb tb gj xn tm pu uh